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Stable-Diffusion-XL-Base-1.0

How to run

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Model Description

Stable Diffusion XL Base 1.0 (SDXL 1.0) is a foundation text-to-image model released by Stability AI.
It is the flagship successor to Stable Diffusion 2.1, designed for photorealism, artistic flexibility, and high-resolution generation.

SDXL 1.0 is a latent diffusion model trained on a broad dataset of images and captions. Compared to prior versions, it improves prompt alignment, visual coherence, and output quality, especially in complex scenes and detailed compositions.

Features

  • High fidelity image generation: sharper details and improved realism.
  • Flexible style range: from photorealistic renders to artistic illustration.
  • Better prompt alignment: improved understanding of nuanced or multi-concept prompts.
  • High resolution support: natively trained for 1024ร—1024 images.
  • Compositional strength: more accurate handling of multiple subjects and fine object placement.

Use Cases

  • Creative content generation (illustrations, art, concept design)
  • Product mockups and marketing visuals
  • Character and environment ideation
  • Storyboarding and visual storytelling
  • Research in generative imaging

Inputs and Outputs

Input:

  • Text prompts (descriptions, concepts, artistic directions)
  • Optional negative prompts to avoid undesired elements

Output:

  • Generated image(s) matching the prompt
  • Default resolution: 1024ร—1024 pixels

License

References

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