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Stable-Diffusion-XL-Base-1.0
How to run
Visit sdk.nexa.ai/model
Model Description
Stable Diffusion XL Base 1.0 (SDXL 1.0) is a foundation text-to-image model released by Stability AI.
It is the flagship successor to Stable Diffusion 2.1, designed for photorealism, artistic flexibility, and high-resolution generation.
SDXL 1.0 is a latent diffusion model trained on a broad dataset of images and captions. Compared to prior versions, it improves prompt alignment, visual coherence, and output quality, especially in complex scenes and detailed compositions.
Features
- High fidelity image generation: sharper details and improved realism.
- Flexible style range: from photorealistic renders to artistic illustration.
- Better prompt alignment: improved understanding of nuanced or multi-concept prompts.
- High resolution support: natively trained for 1024ร1024 images.
- Compositional strength: more accurate handling of multiple subjects and fine object placement.
Use Cases
- Creative content generation (illustrations, art, concept design)
- Product mockups and marketing visuals
- Character and environment ideation
- Storyboarding and visual storytelling
- Research in generative imaging
Inputs and Outputs
Input:
- Text prompts (descriptions, concepts, artistic directions)
- Optional negative prompts to avoid undesired elements
Output:
- Generated image(s) matching the prompt
- Default resolution: 1024ร1024 pixels
License
- Licensed under: CreativeML Open RAIL++-M License
References
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